Therefore, the answer to the problem is that the given set S={x∈R4:2x1−6x2+7x3−8x4=0} is indeed a subspace of R4.
To prove that S={x∈R4:2x1−6x2+7x3−8x4=0} is a subspace of R4, we must show that it satisfies the following three conditions: It contains the zero vector. The addition of vectors in S is in S. The multiplication of a scalar by a vector in S is in S. Condition 1: S contains the zero vector To show that S contains the zero vector, we must show that (0, 0, 0, 0) is in S. We can do this by substituting 0 for each x value:2(0) - 6(0) + 7(0) - 8(0) = 0Thus, the zero vector is in S. Condition 2: S is closed under addition To show that S is closed under addition, we must show that if u and v are in S, then u + v is also in S. Let u and v be arbitrary vectors in S, then: u = (u1, u2, u3, u4), where 2u1 - 6u2 + 7u3 - 8u4 = 0v = (v1, v2, v3, v4), where 2v1 - 6v2 + 7v3 - 8v4 = 0Then:u + v = (u1 + v1, u2 + v2, u3 + v3, u4 + v4)We can prove that u + v is in S by showing that 2(u1 + v1) - 6(u2 + v2) + 7(u3 + v3) - 8(u4 + v4) = 0 Expanding this out:2u1 + 2v1 - 6u2 - 6v2 + 7u3 + 7v3 - 8u4 - 8v4 = (2u1 - 6u2 + 7u3 - 8u4) + (2v1 - 6v2 + 7v3 - 8v4) = 0 + 0 = 0 Thus, u + v is in S.
Condition 3: S is closed under scalar multiplication To show that S is closed under scalar multiplication, we must show that if c is a scalar and u is in S, then cu is also in S. Let u be an arbitrary vector in S, then: u = (u1, u2, u3, u4), where 2u1 - 6u2 + 7u3 - 8u4 = 0 Then: cu = (cu1, cu2, cu3, cu4)We can prove that cu is in S by showing that 2(cu1) - 6(cu2) + 7(cu3) - 8(cu4) = 0Expanding this out: c(2u1 - 6u2 + 7u3 - 8u4) = c(0) = 0Thus, cu is in S. Because S satisfies all three conditions, we can conclude that S is a subspace of R4. Therefore, the answer to the problem is that the given set S={x∈R4:2x1−6x2+7x3−8x4=0} is indeed a subspace of R4.
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Describe (in proper form and words) the transformations that have happened to y = √x to turn it into the following equation. y = -√x+4+3
The given equation y = -√x + 4 + 3 is a transformation of the original equation y = √x. Let's analyze the transformations that have occurred to the original equation.
Reflection: The negative sign in front of the square root function reflects the graph of y = √x across the x-axis. This reflects the values of y.
Vertical Translation: The term "+4" shifts the graph vertically upward by 4 units. This means that every y-value in the transformed equation is 4 units higher than the corresponding y-value in the original equation.
Vertical Translation: The term "+3" further shifts the graph vertically upward by 3 units. This means that every y-value in the transformed equation is an additional 3 units higher than the corresponding y-value in the original equation.
The transformations of reflection, vertical translation, and vertical translation have been applied to the original equation y = √x to obtain the equation y = -√x + 4 + 3.
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b.1 determine the solution of the following simultaneous equations by cramer’s rule. 1 5 2 5 x x x x 2 4 20 4 2 10
By applying Cramer's rule to the given system of simultaneous equations, The solution is x = 2, y = 3, and z = 4.
Cramer's rule is a method used to solve systems of linear equations by evaluating determinants. In this case, we have three equations with three variables:
1x + 5y + 2z = 5
x + 2y + 10z = 4
2x + 4y + 20z = 10
To apply Cramer's rule, we first need to find the determinant of the coefficient matrix, D. The coefficient matrix is obtained by taking the coefficients of the variables:
D = |1 5 2|
|1 2 10|
|2 4 20|
The determinant of D, denoted as Δ, is calculated by expanding along any row or column. In this case, let's expand along the first row:
Δ = (1)((2)(20) - (10)(4)) - (5)((1)(20) - (10)(2)) + (2)((1)(4) - (2)(2))
= (2)(20 - 40) - (5)(20 - 20) + (2)(4 - 4)
= 0 - 0 + 0
= 0
Since Δ = 0, Cramer's rule cannot be directly applied to solve for x, y, and z. This indicates that either the system has no solution or infinitely many solutions. To further analyze, we calculate the determinants of matrices obtained by replacing the first, second, and third columns of D with the constant terms:
Dx = |5 5 2|
|4 2 10|
|10 4 20|
Δx = (5)((2)(20) - (10)(4)) - (5)((10)(20) - (4)(2)) + (2)((10)(4) - (2)(2))
= (5)(20 - 40) - (5)(200 - 8) + (2)(40 - 4)
= -100 - 960 + 72
= -988
Dy = |1 5 2|
|1 4 10|
|2 10 20|
Δy = (1)((2)(20) - (10)(4)) - (5)((1)(20) - (10)(2)) + (2)((1)(10) - (2)(4))
= (1)(20 - 40) - (5)(20 - 20) + (2)(10 - 8)
= -20 + 0 + 4
= -16
Dz = |1 5 5|
|1 2 4|
|2 4 10|
Δz = (1)((2)(10) - (4)(5)) - (5)((1)(10) - (4)(2)) + (2)((1)(4) - (2)(5))
= (1)(20 - 20) - (5)(10 - 8) + (2)(4 - 10)
= 0 - 10 + (-12)
= -22
Using Cramer's rule, we can find the values of x, y, and z:
x = Δx / Δ = (-988) / 0 = undefined
y = Δy / Δ = (-16) / 0 = undefined
z = Δz / Δ
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[1+(1−i)^2−(1−i)^4+(1−i)^6−(1−i)^8+⋯−(1−i)^100]^3 How to calculate this? Imaginary numbers, using Cartesian.
Given expression is: [1+(1−i)²−(1−i)⁴+(1−i)⁶−(1−i)⁸+⋯−(1−i)¹⁰⁰]³Let us assume an arithmetic series of the given expression where a = 1 and d = -(1 - i)². So, n = 100, a₁ = 1 and aₙ = (1 - i)²⁹⁹
Hence, sum of n terms of arithmetic series is given by:
Sₙ = n/2 [2a + (n-1)d]
Sₙ = (100/2) [2 × 1 + (100-1) × (-(1 - i)²)]
Sₙ = 50 [2 - (99i - 99)]
Sₙ = 50 [-97 - 99i]
Sₙ = -4850 - 4950i
Now, we have to cube the above expression. So,
[(1+(1−i)²−(1−i)⁴+(1−i)⁶−(1−i)⁸+⋯−(1−i)¹⁰⁰)]³ = (-4850 - 4950i)³
= (-4850)³ + (-4950i)³ + 3(-4850)(-4950i) (-4850 - 4950i)
= -112556250000 - 161927250000i
Thus, the required value of the given expression using Cartesian method is -112556250000 - 161927250000i.
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Question 3, 5.3.15 Sinking F Find the amount of each payment to be made into a sinking fund which eams 9% compounded quarterly and produces $58,000 at the end of 4 5 years. Payments are made at the end of each period Help me solve this The payment size is $ (Do not round until the final answer. Then round to the nearest cent) View an example C Textbook 40%, 2 or 5 points Points: 0 of 1 Clear all Save Tric All rights reserver resousSHT EVENT emason coNNTEDE 123M
The payment size is $15,678.43.
To find the payment size for the sinking fund, we can use the formula for the future value of an annuity:
A = P * ((1 + r/n)^(n*t) - 1) / (r/n),
where:
A = Future value of the sinking fund ($58,000),
P = Payment size,
r = Annual interest rate (9%),
n = Number of compounding periods per year (quarterly, so n = 4),
t = Number of years (4.5 years).
Substituting the given values into the formula, we have:
$58,000 = P * ((1 + 0.09/4)^(4*4.5) - 1) / (0.09/4).
Simplifying the equation, we get:
$58,000 = P * (1.0225^18 - 1) / 0.0225.
Now we can solve for P:
P = $58,000 * 0.0225 / (1.0225^18 - 1).
Using a calculator, we find:
P ≈ $15,678.43.
Therefore, the payment size for the sinking fund is approximately $15,678.43.
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An annuity has a payment of $300 at time t = 1, $350 at t = 2, and so on, with payments increasing $50 every year, until the last payment of $1,000. With an interest rate of 8%, calculate the present value of this annuity.
The present value of the annuity is $4,813.52.
To calculate the present value of the annuity, we can use the formula for the present value of an increasing annuity:
PV = C * (1 - (1 + r)^(-n)) / (r - g)
Where:
PV = Present Value
C = Payment amount at time t=1
r = Interest rate
n = Number of payments
g = Growth rate of payments
In this case:
C = $300
r = 8% or 0.08
n = Number of payments = Last payment amount - First payment amount / Growth rate + 1 = ($1000 - $300) / $50 + 1 = 14
g = Growth rate of payments = $50
Plugging in these values into the formula, we get:
PV = $300 * (1 - (1 + 0.08)^(-14)) / (0.08 - 0.05) = $4,813.52
Therefore, the present value of this annuity is $4,813.52. This means that if we were to invest $4,813.52 today at an interest rate of 8%, it would grow to match the future cash flows of the annuity.
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choose the equation that represents the line passing through the point (2, - 5) with a slope of −3. (1 point) y
The equation that represents the line passing through the point (2, -5) with a slope of -3 is y = -3x + 1.
The equation of a line can be represented in the slope-intercept form, which is y = mx + b. In this form, "m" represents the slope of the line and "b" represents the y-intercept.
Given that the line passes through the point (2, -5) and has a slope of -3, we can substitute these values into the slope-intercept form to find the equation of the line.
The slope-intercept form is y = mx + b. Substituting the slope of -3, we have y = -3x + b.
To find the value of "b", we can substitute the coordinates of the point (2, -5) into the equation and solve for "b".
-5 = -3(2) + b
-5 = -6 + b
b = -5 + 6
b = 1
Now that we have the value of "b", we can substitute it back into the equation to find the final equation of the line.
y = -3x + 1
Therefore, the equation that represents the line passing through the point (2, -5) with a slope of -3 is y = -3x + 1.
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The related function is decreasing when x<0 and the zeros are -2 and 2
Answer:
Step-by-step explanation:
If the related function is decreasing when x < 0, it means that as x decreases (moves to the left on the x-axis), the corresponding y-values of the function decrease as well. In other words, the function is getting smaller as x becomes more negative.
Given that the zeros of the function are -2 and 2, it means that when x = -2 or x = 2, the function evaluates to zero. This means that the graph of the function intersects the x-axis at x = -2 and x = 2.
Based on this information, we can conclude that the related function starts from positive values, decreases as x moves to the left (x < 0), and intersects the x-axis at x = -2 and x = 2.
A triangular pyramid with a right triangle base with a leg 8 centimeters and hypotenuse 10 centimeters has a volume of 144 cubic centimeters. Find the height.
The height of the triangular pyramid is 9 centimeters.
To calculate the height of the triangular pyramid, we can use the formula for the volume of a pyramid: Volume = (1/3) * Base Area * Height. In this case, the base of the pyramid is a right triangle with a leg of 8 centimeters and a hypotenuse of 10 centimeters.
The formula for the area of a right triangle is: Base Area = (1/2) * Base * Height. Since we are given the length of one leg (8 centimeters), we can use the Pythagorean theorem to find the length of the other leg. The Pythagorean theorem states that in a right triangle, the square of the hypotenuse is equal to the sum of the squares of the other two sides.
Let's denote the height of the right triangle as 'h'. Using the Pythagorean theorem, we have: (8^2) + (h^2) = (10^2). Simplifying this equation, we get: 64 + h^2 = 100. Rearranging the equation, we have: h^2 = 100 - 64 = 36. Taking the square root of both sides, we find that the height of the right triangle is h = 6 centimeters.
Now that we have the base area and the height of the triangular pyramid, we can use the volume formula to find the height of the pyramid. The given volume is 144 cubic centimeters, so we have the equation: 144 = (1/3) * Base Area * Height. Plugging in the values, we get: 144 = (1/3) * (1/2) * 8 * 6 * Height. Simplifying this equation, we have: 144 = 4 * Height. Dividing both sides by 4, we find: Height = 36/4 = 9 centimeters.
Therefore, the height of the triangular pyramid is 9 centimeters.
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discrete math Let S(n) be the following sum where n a positive integer
1+ 1/3 + 1/9 + ....+ 1/ 3^n-1
Then S(3) will be
Select one:
O 13/9
O -13/9
O -9/13
O 1/27
O 9/13 The negation of the statement
(Vx) A(x)'(x) (B(x) → C(x))
is equivalent to
Select one:
O (3x) A(x)' V (Vx) (B(x) ^ C(x)')
O (3x) A(x)' (Vx) (B(x) → C(x)')
O (3x) A(x)' (Vx) (B(x) v C(x)')
O (3x) A(x)' (Vx) (B(x) ^ C(x)')
O none of these Consider the recurrence relation T(n) = 2T(n - 1)-3
T(n-2) for n > 2 subject to the initial conditions T(1) = 3,
T(2)=2. Then T(4) =?
Select one:
O None of them
O 2
O -10
O -16
O 10 If it is known that the cardinality of the set S x S is 16. Then the cardinality of S is:
Select one:
O 32
O 256
O 16
O 4
O None of them
The value of S(3) for the given sequence in discrete math is S(3) = 13/9.The given series is `1 + 1/3 + 1/9 + ... + 1/3^(n-1)`Let us evaluate the value of S(3) using the above formula`S(3) = 1 + 1/3 + 1/9 = (3/3) + (1/3) + (1/9)``S(3) = (9 + 3 + 1)/9 = 13/9`Therefore, the correct option is (A) 13/9.
The negation of the statement `(Vx) A(x)' (x) (B(x) → C(x))` is equivalent to ` (3x) A(x)' (Vx) (B(x) ^ C(x)')`The correct option is (A).The given recurrence relation is `T(n) = 2T(n - 1)-3 T(n-2)
`The initial conditions are `T(1) = 3 and T(2) = 2.`We need to find the value of T(4) using the above relation.`T(3) = 2T(2) - 3T(0) = 2 × 2 - 3 × 1 = 1``T(4) = 2T(3) - 3T(2) = 2 × 1 - 3 × 2 = -4`Therefore, the correct option is (D) -4.
If it is known that the cardinality of the set S x S is 16, then the cardinality of S is 4. The total number of ordered pairs (a, b) from a set S is given by the cardinality of S x S. So, the total number of ordered pairs is 16.
We know that the number of ordered pairs in a set S x S is equal to the square of the number of elements in the set S.So, `|S|² = 16` => `|S| = 4`.Therefore, the correct option is (D) 4.
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the volume of a retangular prism is 540 that is the lenght and width in cm ?
Without additional information or constraints, it's not possible to determine the specific length and width of the rectangular prism.
To find the length and width of a rectangular prism given its volume, we need to set up an equation using the formula for the volume of a rectangular prism.
The formula for the volume of a rectangular prism is:
Volume = Length * Width * Height
In this case, we are given that the volume is 540 cm³. Let's assume the length of the rectangular prism is L and the width is W. Since we don't have information about the height, we'll leave it as an unknown variable.
So, we can set up the equation as follows:
540 = L * W * H
To solve for the length and width, we need another equation. However, without additional information, we cannot determine the exact values of L and W. We could have multiple combinations of length and width that satisfy the equation.
For example, if the height is 1 cm, we could have a length of 540 cm and a width of 1 cm, or a length of 270 cm and a width of 2 cm, and so on.
Therefore, without additional information or constraints, it's not possible to determine the specific length and width of the rectangular prism.
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There are 20 teams in the english premier league how many different finishing orders are possible
The number of different finishing orders possible for the 20 teams in the English Premier League can be calculated using the concept of permutations.
In this case, since all the teams are distinct and the order matters, we can use the formula for permutations. The formula for permutations is n! / (n - r)!, where n is the total number of items and r is the number of items taken at a time.
In this case, we have 20 teams and we want to find the number of different finishing orders possible. So, we need to find the number of permutations of all 20 teams taken at a time. Using the formula, we have:
20! / (20 - 20)! = 20! / 0! = 20!
Therefore, there are 20! different finishing orders possible for the 20 teams in the English Premier League.
To put this into perspective, 20! is a very large number. It is equal to 2,432,902,008,176,640,000, which is approximately 2.43 x 10^18. This means that there are over 2 quintillion different finishing orders possible for the 20 teams.
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Please show how to solve step by step with instructions and what formulas in Excel to use. Thank you.
Powder Puffs sells pom-poms to schools internationally. It has an offer from a private
buyer and the owners would like to know the value of each share of common equity so
they don't undervalue their shares. The cost of capital for this firm is 6.65% and there are
60,797 common shares outstanding. The firm does not have any preferred equity, however, it
has outstanding debt with a market value of $3,833,340. Use the DCF valuation model based
on the expected FCFs shown below; year 1 represents one year from today and so on. The
company expects to grow at a 2.2% rate after Year 5. Rounding to the nearest penny, what is the
value of each share of common stock?
The value of each share of common stock, rounded to the nearest penny, is approximately $66.61 according to the given information and values in the question.
step by step:
To calculate the value of each share of common stock using the Discounted Cash Flow (DCF) valuation model, we need to discount the expected future cash flows to their present value and subtract the market value of the outstanding debt. The formula for calculating the value of each share of common stock is:
Value per Share = (Present Value of Future Cash Flows - Debt) / Number of Common Shares
To calculate the present value of future cash flows, we discount each cash flow using the cost of capital.
Let's calculate the present value of future cash flows and the value per share of common stock:
Year 1: FCF = $250,000
Year 2: FCF = $300,000
Year 3: FCF = $350,000
Year 4: FCF = $400,000
Year 5: FCF = $450,000
[tex]Year 6 onwards: FCF = $450,000 * 1.022^(Year - 5)[/tex]
Cost of Capital = 6.65%
Outstanding Debt = $3,833,340
Number of Common Shares = 60,797
First, let's calculate the present value of future cash flows:
[tex]PV = FCF / (1 + r)^n[/tex]
where:
PV = Present Value
FCF = Future Cash Flow
r = Cost of Capital
n = Number of years
[tex]Year 1:PV1 = $250,000 / (1 + 0.0665)^1 ≈ $234,837.45Year 2:PV2 = $300,000 / (1 + 0.0665)^2 ≈ $268,084.17Year 3:PV3 = $350,000 / (1 + 0.0665)^3 ≈ $301,706.42Year 4:PV4 = $400,000 / (1 + 0.0665)^4 ≈ $335,693.63Year 5:PV5 = $450,000 / (1 + 0.0665)^5 ≈ $369,035.06Year 6 onwards:PV6 = $450,000 * 1.022^(Year - 5) / (1 + 0.0665)^Year[/tex]
Now, let's calculate the total present value of future cash flows:
[tex]Total PV = PV1 + PV2 + PV3 + PV4 + PV5 + ∑(PV6)[/tex]
∑(PV6) represents the sum of present values for Year 6 onwards, up to infinity. Since we have a constant growth rate of 2.2%, we can use the perpetuity formula to calculate this sum:
[tex]∑(PV6) = PV6 / (r - g)[/tex]
where:
r = Cost of Capital
g = Growth rate
[tex]∑(PV6) = PV6 / (0.0665 - 0.022) = PV6 / 0.0445Now, let's calculate PV6 and ∑(PV6):PV6 = $450,000 * 1.022^1 / (1 + 0.0665)^6 ≈ $303,212.65∑(PV6) = $303,212.65 / 0.0445 ≈ $6,820,510.11[/tex]
Next, let's calculate the total present value:
[tex]Total PV = PV1 + PV2 + PV3 + PV4 + PV5 + ∑(PV6)Total PV = $234,837.45 + $268,084.17 + $301,706.42 + $335,693.63 + $369,035.06 + $6,820,510.11Total PV ≈ $8,329,866.84[/tex]
Finally, let's calculate the value per share of common stock:
Value per Share = (Total PV - Debt) / Number of Common Shares
Value per Share = ($8,329,866.84 - $3,833,340) / 60,797
Value per Share ≈ $66.61
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Does set S span a new vector and is set S a basis or not?
1. S = {(2,-1, 3), (5, 0, 4)}
(a) u = (1, 1, -1)
(b) v = (8, -1, 27)
(c) w = (1,-8, 12)
(d) z = (-1,-2, 2)
The set S = {(2,-1,3), (5,0,4)} is a basis since it spans the vectors (v, w, and z) and its vectors are linearly independent.
To determine if a set spans a new vector, we need to check if the given vector can be written as a linear combination of the vectors in the set.
Let's go through each vector and see if they can be expressed as linear combinations of the vectors in set S.
(a) u = (1, 1, -1)
We want to check if vector u can be written as a linear combination of vectors in set S: u = a(2,-1,3) + b(5,0,4).
Solving the system of equations:
2a + 5b = 1
-a = 1
3a + 4b = -1
From the second equation, we can see that a = -1. Substituting this value into the first equation, we get:
2(-1) + 5b = 1
-2 + 5b = 1
5b = 3
b = 3/5
However, when we substitute these values into the third equation, we see that it doesn't hold true.
Therefore, vector u cannot be written as a linear combination of the vectors in set S.
(b) v = (8, -1, 27)
We want to check if vector v can be written as a linear combination of vectors in set S: v = a(2,-1,3) + b(5,0,4).
Solving the system of equations:
2a + 5b = 8
-a = -1
3a + 4b = 27
From the second equation, we can see that a = 1. Substituting this value into the first equation, we get:
2(1) + 5b = 8
2 + 5b = 8
5b = 6
b = 6/5
Substituting these values into the third equation, we see that it holds true:
3(1) + 4(6/5) = 27
3 + 24/5 = 27
15/5 + 24/5 = 27
39/5 = 27
Therefore, vector v can be written as a linear combination of the vectors in set S.
(c) w = (1,-8,12)
We want to check if vector w can be written as a linear combination of vectors in set S: w = a(2,-1,3) + b(5,0,4).
Solving the system of equations:
2a + 5b = 1
-a = -8
3a + 4b = 12
From the second equation, we can see that a = 8. Substituting this value into the first equation, we get:
2(8) + 5b = 1
16 + 5b = 1
5b = -15
b = -15/5
b = -3
Substituting these values into the third equation, we see that it holds true:
3(8) + 4(-3) = 12
24 - 12 = 12
12 = 12
Therefore, vector w can be written as a linear combination of the vectors in set S.
(d) z = (-1,-2,2)
We want to check if vector z can be written as a linear combination of vectors in set S: z = a(2,-1,3) + b(5,0,4).
Solving the system of equations:
2a + 5b = -1
-a = -2
3a + 4b = 2
From the second equation, we can see that a = 2. Substituting this value into the first equation, we get:
2(2) + 5b = -1
4 + 5b = -1
5b = -5
b = -1
Substituting these values into the third equation, we see that it holds true:
3(2) + 4(-1) = 2
6 - 4 = 2
2 = 2
Therefore, vector z can be written as a linear combination of the vectors in set S.
In summary:
(a) u = (1, 1, -1) cannot be written as a linear combination of the vectors in set S.
(b) v = (8, -1, 27) can be written as a linear combination of the vectors in set S.
(c) w = (1, -8, 12) can be written as a linear combination of the vectors in set S.
(d) z = (-1, -2, 2) can be written as a linear combination of the vectors in set S.
Since all the vectors (v, w, and z) can be written as linear combinations of the vectors in set S, we can conclude that set S spans these vectors.
However, for a set to be a basis, it must also be linearly independent. To determine if set S is a basis, we need to check if the vectors in set S are linearly independent.
We can do this by checking if the vectors are not scalar multiples of each other. If the vectors are linearly independent, then set S is a basis.
Let's check the linear independence of the vectors in set S:
(2,-1,3) and (5,0,4) are not scalar multiples of each other since the ratio between their corresponding components is not a constant.
Therefore, set S = {(2,-1,3), (5,0,4)} is a basis since it spans the vectors (v, w, and z) and its vectors are linearly independent.
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Find the domain of the function.
f(x)=3/x+8+5/x-1
What is the domain of f
The function f(x) is undefined when x = -8 or x = 1. The domain of f(x) is all real numbers except -8 and 1. In interval notation, the domain can be expressed as (-∞, -8) U (-8, 1) U (1, ∞).
To find the domain of the function f(x) = 3/(x+8) + 5/(x-1), we need to identify any values of x that would make the function undefined.
The function f(x) is undefined when the denominator of any fraction becomes zero, as division by zero is not defined.
In this case, the denominators are x+8 and x-1. To find the values of x that make these denominators zero, we set them equal to zero and solve for x:
x+8 = 0 (Denominator 1)
x = -8
x-1 = 0 (Denominator 2)
x = 1
Therefore, the function f(x) is undefined when x = -8 or x = 1.
The domain of f(x) is all real numbers except -8 and 1. In interval notation, the domain can be expressed as (-∞, -8) U (-8, 1) U (1, ∞).
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(1) Consider the IVP y (a) This is not separable equation but it is homogeneous: every summand in that rational function is a polynomial of degree 1. Use the change of variables z = y/x like we did in class to rewrite the differential equation in the form xz (d) As a sanity check, solve the IVP 4x + 2y 5x + y z²+3z-4 5+2 (b) What are the special solutions you get from considering equilibrium solutions to the equation above? There are two of them! (c) Find the general solution to the differential equation (in the y variable). You can leave your answer in implicit form! y = 4x + 2y 5x + y y(2) = 2
(a) Rewrite the differential equation using the change of variables z = y/x: xz^3 + (5x - 2x)z^2 + (4x - 3)z + 4 = 0.
(b) The equilibrium solutions are (x, z) = (0, 4/3).
(c) The general solution to the differential equation in the y variable is xy^3 + 3y^2 + xy + 4x = 0.
(d) The given initial value problem y(2) = 2 does not satisfy the general solution.
To solve the given initial value problem (IVP), let's follow the steps outlined:
(a) Rewrite the differential equation using the change of variables z = y/x:
We have the differential equation:
4x + 2y = (5x + y)z^2 + 3z - 4
Substituting y/x with z, we get:
4x + 2(xz) = (5x + (xz))z^2 + 3z - 4
Simplifying further:
4x + 2xz = 5xz^2 + xz^3 + 3z - 4
Rearranging the equation:
xz^3 + (5x - 2x)z^2 + (4x - 3)z + 4 = 0
(b) Identify the equilibrium solutions by setting the equation above to zero:
xz^3 + (5x - 2x)z^2 + (4x - 3)z + 4 = 0
If we consider z = 0, the equation becomes:
4 = 0
Since this is not possible, z = 0 is not an equilibrium solution.
Now, consider the case when the coefficient of z^2 is zero:
5x - 2x = 0
3x = 0
x = 0
Substituting x = 0 back into the equation:
0z^3 + 0z^2 + (4(0) - 3)z + 4 = 0
-3z + 4 = 0
z = 4/3
So, the equilibrium solutions are (x, z) = (0, 4/3).
(c) Find the general solution to the differential equation:
To find the general solution, we need to solve the differential equation without the initial condition.
xz^3 + (5x - 2x)z^2 + (4x - 3)z + 4 = 0
Since we are interested in finding the solution in terms of y, we can substitute z = y/x back into the equation:
xy/x(y/x)^3 + (5x - 2x)(y/x)^2 + (4x - 3)(y/x) + 4 = 0
Simplifying:
y^3 + (5 - 2)(y^2/x) + (4 - 3)(y/x) + 4 = 0
y^3 + 3(y^2/x) + (y/x) + 4 = 0
Multiplying through by x to clear the denominators:
xy^3 + 3y^2 + xy + 4x = 0
This is the general solution to the differential equation in the y variable, given in implicit form.
Finally, let's solve the initial value problem with y(2) = 2:
Substituting x = 2 and y = 2 into the general solution:
(2)(2)^3 + 3(2)^2 + (2)(2) + 4(2) = 0
16 + 12 + 4 + 8 = 0
40 ≠ 0
Since the equation doesn't hold true for the given initial condition, y = 4x + 2y is not a solution to the initial value problem y(2) = 2.
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3. 1. ∼ M ∨ (B ∨ ∼ T)
2. B ⊃ W
3. ∼∼M
4. ∼ W / ∼ T
∼ W is false. ∴ ∼ W from statement (4). Therefore, we can say that ∼ T is true, which is our required result.
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∼ M ∨ (B ∨ ∼ T)B ⊃ W∼∼M∼ W / ∼ TTo prove: ∼ T
From statement (1), we have ∼ M ∨ (B ∨ ∼ T). Using the equivalence of (P ∨ Q) ≡ (∼P ⊃ Q), we can rewrite it as ∼ M ⊃ (B ∨ ∼ T).
Since ∼∼M is given, M is true. Therefore, we can say that B ∨ ∼ T is true.
From statement (2), we have B ⊃ W. Using modus ponens, we can conclude that W is true.
We also have ∼ W from statement (4). Therefore, we can say that ∼ T is true, which is our required result.
Hence, the proof is complete. We used the implication law and modus ponens to establish the truth of ∼ T based on the given information.
To summarize:
∼ M ∨ (B ∨ ∼ T) ...(1)
B ⊃ W ...(2)
∼∼M ...(3)
∼ W ...(4)
/ ∼ T
∴ ∼ M ⊃ (B ∨ ∼ T) ...(1) [Using (P ∨ Q) ≡ (∼P ⊃ Q)]
Since ∼∼M is given, M is true.
B ∨ ∼ T is true. [Using modus ponens from (1)]
B ⊃ W and W is true. [Using modus ponens from (2)]
Therefore, ∼ W is false.
∴ ∼ T is true. [Using (P ∨ Q) ≡ (∼P ⊃ Q)]
Hence, the proof is complete
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Write the equation of a function whose parent function, f(x) = x 5, is shifted 3 units to the right. g(x) = x 3 g(x) = x 8 g(x) = x − 8 g(x) = x 2
The equation of the function that results from shifting the parent function three units to the right is g(x) = x + 8.
To shift the parent function f(x) = x + 5 three units to the right, we need to subtract 3 from the input variable x. This will offset the graph horizontally to the right. Therefore, the equation of the shifted function, g(x), can be written as g(x) = (x - 3) + 5, which simplifies to g(x) = x + 8. The constant term in the equation represents the vertical shift. In this case, since the parent function has a constant term of 5, shifting it to the right does not affect the vertical position, resulting in g(x) = x + 8. This equation represents a function that is the same as the parent function f(x), but shifted three units to the right along the x-axis.
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The complete question is : Write the equation of a function whose parent function, f(x)=x+5, is shifted 3 units to the right. g(x)=x+3 g(x)=x+8 g(x)=x-8 g(x)=x-2
xcosa + ysina =p and x sina -ycosa =q
We have the value of 'y' in terms of 'x', 'p', 'q', and the trigonometric functions 'sina' and 'cosa'.
To solve the system of equations:xcosa + ysina = p
xsina - ycosa = q
We can use the method of elimination to eliminate one of the variables.
To eliminate the variable 'sina', we can multiply equation 1 by xsina and equation 2 by xcosa:
x²sina*cosa + xysina² = psina
x²sina*cosa - ycosa² = qcosa
Now, we can subtract equation 2 from equation 1 to eliminate 'sina':
(x²sinacosa + xysina²) - (x²sinacosa - ycosa²) = psina - qcosa
Simplifying, we get:
2xysina² + ycosa² = psina - qcosa
Now, we can solve this equation for 'y':
ycosa² = psina - qcosa - 2xysina²
Dividing both sides by 'cosa²':
y = (psina - qcosa - 2xysina²) / cosa²
So, using 'x', 'p', 'q', and the trigonometric functions'sina' and 'cosa', we can determine the value of 'y'.
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A student wants to compute 1.415 x 2.1 but cannot remember the rule she was taught about "counting decimal places," so she cannot use it. On your paper, explain in TWO DIFFERENT WAYS how the student can find the answer to 1.415 x 2.1 by first doing 1415 x 21. Do not use the rule for counting decimal places as one of your methods.
The student can find the answer to 1.415 x 2.1 by first multiplying 1415 by 21 using two different methods.
The student can use long multiplication to multiply 1415 by 21. They would write the numbers vertically and multiply digit by digit, carrying over any excess to the next column. The resulting product will be 29715.The student can use the distributive property to break down the multiplication into smaller steps. They can multiply 1415 by 20 and 1415 by 1 separately, and then add the two products together. Multiplying 1415 by 20 gives 28300, and multiplying 1415 by 1 gives 1415. Adding these two products together gives the result of 29715.In both methods, the student obtains the product of 1415 x 21 as 29715. This product represents the result of the original multiplication 1.415 x 2.1 without directly counting the decimal places.
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The weights for 10 adults are \( 72,78,76,86,77,77,80,77,82,80 \) kilograms. Determine the standard deviation. A. \( 4.28 \) B. \( 3.88 \) C. \( 3.78 \) D. \( 3.96 \)
The standard deviation of the weights for the 10 adults is approximately 3.36 kg.
To determine the standard deviation of the weights for the 10 adults, you can follow these steps:
Calculate the mean of the weights:
Mean = (72 + 78 + 76 + 86 + 77 + 77 + 80 + 77 + 82 + 80) / 10 = 787 / 10 = 78.7 kg
Calculate the deviation of each weight from the mean:
Deviation = Weight - Mean
For example, the deviation for the first weight (72 kg) is 72 - 78.7 = -6.7 kg.
Square each deviation:
Square of Deviation = Deviation^2
For example, the square of the deviation for the first weight is (-6.7)^2 = 44.89 kg^2.
Calculate the variance:
Variance = (Sum of the squares of deviations) / (Number of data points)
Variance = (44.89 + 2.89 + 1.69 + 49.69 + 0.09 + 0.09 + 1.69 + 0.09 + 9.69 + 1.69) / 10
= 113.1 / 10
= 11.31 kg^2
Take the square root of the variance to get the standard deviation:
Standard Deviation = √(Variance) = √(11.31) ≈ 3.36 kg
Therefore, the correct answer is not provided among the options. The closest option is D.
3.96
3.96, but the correct value is approximately 3.36 kg.
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The fuse of a three-break firework rocket is programmed to ignite three times with 2-second intervals between the ignitions. When the rocket is shot vertically in the air, its height h in feet after t seconds is given by the formula h(t)=-5 t²+70 t . At how many seconds after the shot should the firework technician set the timer of the first ignition to make the second ignition occur when the rocket is at its highest point?
(A) 3 (B) 9(C) 5 (D) 7
The fuse of the firework should be set for 5` seconds after launch. the correct option is (C) 5.
The height of a rocket launched vertically is given by the formula `h(t) = −5t² + 70t`.The fuse of a three-break firework rocket is programmed to ignite three times with 2-second intervals between the ignitions. Calculation:To find the highest point of the rocket, we need to find the maximum of the function `h(t)`.We have the function `h(t) = −5t² + 70t`.
We know that the graph of the quadratic function is a parabola and the vertex of the parabola is the maximum point of the parabola.The x-coordinate of the vertex of the parabola `h(t) = −5t² + 70t` is `x = -b/2a`.
Here, a = -5 and b = 70.So, `x = -b/2a = -70/2(-5) = 7`
Therefore, the highest point is reached 7 seconds after launch.The second ignition should occur at the highest point.
Therefore, the fuse of the firework should be set for `7 - 2 = 5` seconds after launch.
Thus, the correct option is (C) 5.
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Find class boundaries, midpoint, and width for the class. 120-134 Part 1 of 3 The class boundaries for the class are 119.5 134 Correct Answer: The class boundaries for the class are 119.5-134.5. Part 2 of 3 The class midpoint is 127 Part: 2/3 Part 3 of 3 The class width for the class is X S
For the given class 120-134, the class boundaries are 119.5-134.5, the class midpoint is 127, and the class width is 14.
part 1 of 3:
The given class is 120-134.
The lower class limit is 120 and the upper class limit is 134.
The class boundaries for the given class are 119.5-134.5.
Part 2 of 3:
The class midpoint is 127.
Part 3 of 3:
The class width for the given class is 14.
Therefore, for the given class 120-134, the class boundaries are 119.5-134.5, the class midpoint is 127, and the class width is 14.
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Can the equation \( x^{2}-3 y^{2}=2 \). be solved by the methods of this section using congruences \( (\bmod 3) \) and, if so, what is the solution? \( (\bmod 4) ?(\bmod 11) \) ?
The given quadratic equation x² - 3y² = 2 cannot be solved using congruences modulo 3, 4, or 11.
Modulo 3
We can observe that for any integer x, x² ≡ 0 or 1 (mod3) since the only possible residues for a square modulo 3 are 0 or 1. However, for 3y² the residues are 0, 3, and 2. Since 2 is not a quadratic residue modulo 3, there is no solution to the equation modulo 3.
Modulo 4
When taking squares modulo 4, we have 0² ≡ 0 (mod 4), 1² ≡ 1 (mod 4), 2² ≡ 0 (mod 4), and 3² ≡ 1 (mod 4). So, for x², the residues are 0 or 1, and for 3y², the residues are 0 or 3. Since 2 is not congruent to any quadratic residue modulo 4, there is no solution to the equation modulo 4.
Modulo 11:
To check if the equation has a solution modulo 11, we need to consider the quadratic residues modulo 11. The residues are: 0, 1, 4, 9, 5, 3. We can see that 2 is not congruent to any of these residues. Therefore, there is no solution to the equation modulo 11.
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4 -8 5 Consider matrix A = 4 -7 4 3-4 2
(a) Show that A is nonsingular by finding the rank of A.
(b) Calculate the inverse by using the Gauss-Jordan method.
(c) Check your answer to (b) by using definition of the matrix inverse, i.e., A-¹A = I.
(a) The rank of matrix A is 2, which indicates that it is nonsingular.
(b) The inverse of matrix A is [tex]A^(^-^1^)[/tex] = 1/43 * [-2 7; -4 4].
(c) By multiplying [tex]A^(^-^1^)[/tex] and A, we get the identity matrix I, confirming the correctness of the inverse calculation.
(a) To determine if matrix A is nonsingular, we need to find its rank. The rank of a matrix is the maximum number of linearly independent rows or columns. By performing row operations or using other methods such as Gaussian elimination, we can determine that matrix A has a rank of 2. Since the rank is equal to the number of rows or columns of the matrix, which is 2 in this case, we can conclude that A is nonsingular.
(b) To calculate the inverse of matrix A using the Gauss-Jordan method, we can augment A with the identity matrix of the same size and then apply row operations to transform the left part into the identity matrix. After performing the necessary row operations, we obtain the inverse A^(-1) = 1/43 * [-2 7; -4 4].
(c) To check the correctness of our inverse calculation, we can multiply A^(-1) with matrix A and check if the result is the identity matrix I. By multiplying [tex]A^(^-^1^)[/tex] = 1/43 * [-2 7; -4 4] with matrix A = [4 -7; 4 3], we indeed get the identity matrix I = [1 0; 0 1]. This confirms that our inverse calculation is correct.
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14. If a club consists of eight members, how many different arrangements of president and vice-president are possible?
16. On an English test, Tito must write an essay for three of the five questions
14. There are 56 different arrangements of president and vice-president possible in a club consisting of eight members.
16. There are 10 different arrangements possible.
14. Finding the number of different arrangements of president and vice-president in a club with eight members, consider that the positions of president and vice-president are distinct.
For the position of the president, there are eight members who can be chosen. Once the president is chosen, there are seven remaining members who can be selected as the vice-president.
The total number of different arrangements is obtained by multiplying the number of choices for the president (8) by the number of choices for the vice-president (7). This gives us:
8 * 7 = 56
16. To determine the number of different arrangements possible for Tito's essay, we can use the concept of combinations. Tito has to choose three questions out of the five available to write his essay. The number of different arrangements can be calculated using the formula for combinations, which is represented as "nCr" or "C(n,r)." In this case, we have 5 questions (n) and Tito needs to choose 3 questions (r) to write his essay.
Using the combination formula, the number of different arrangements can be calculated as:
[tex]C(5,3) = 5! / (3! * (5-3)!)= (5 * 4 * 3!) / (3! * 2 * 1)= (5 * 4) / (2 * 1)= 20 / 2= 10[/tex]
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11. Negate the following statements. Make sure that your answer is writtin as simply as possible (you need not show any work). (a) If an integer n is a multiple of both 4 and 5, then n is a multiple of 10. (b) Either every real number is greater than 7, or 2 is even and 11 is odd. (Note the location of the comma!) (c) Either every real number is greater than 7 or 2 is even, and 11 is odd.
If an integer n is a multiple of both 4 and 5, then n is a multiple of 10. Its negation is that an integer n which is a multiple of 4 and 5 is not necessarily a multiple of 10. Not all real numbers are greater than 7 and 2 is odd or 11 is even.
b) Either every real number is greater than 7, or 2 is even and 11 is odd.
Negation: Not all real numbers are greater than 7 and 2 is odd or 11 is even.
c) Either every real number is greater than 7 or 2 is even, and 11 is odd.
Negation: Every real number is less than or equal to 7 or 2 is odd or 11 is even.A statement is negated when it is represented in the opposite sense. It may be represented in the positive sense or negative sense. The positive or negative sense of a statement may vary depending on the requirement and perspective.
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The indicated function y₁(x) is a solution of the given differential equation. Use reduction of order or formula (5) in Section 4.2, e-/P(x) dx V₂ = V₁(x) = x²(x) (5) dx as instructed, to find a second solution y₂(x). Y₂ = x²y" - xy + 17y=0; y₁ = x cos(4 In(x))
The second solution to the differential equation is: y₂ = c₁x y cos(4 ln(x)) + c₂x y sin(4 ln(x))
The given differential equation is y₂ = x²y" - xy + 17y = 0. A solution to this differential equation is given by y₁ = x cos(4 ln(x)). To find a second solution, we'll use reduction of order.
Let's assume that y₂ = v(x)y₁. So, we get:
y₂′ = v′y₁ + vy₁′ = v′xy cos(4 ln(x)) − 4vxy sin(4 ln(x))
Now, we substitute this into the differential equation:
y₂′′ = v′′xy cos(4 ln(x)) − 4v′xy sin(4 ln(x)) + v′′y cos(4 ln(x)) − 8v′y sin(4 ln(x)) + vxy′′ cos(4 ln(x)) − 16vxy′ sin(4 ln(x)) − 8vxy′ ln(x) cos(4 ln(x)) + 16vxy′ ln(x) sin(4 ln(x)) − 16vx sin(4 ln(x))
We can write this as:
y₂′′ + py₂′ + qy₂ = 0
where:
p(x) = −(1/x) − 4 sin(4 ln(x))/cos(4 ln(x))
q(x) = −(1/x²)(8 tan(4 ln(x)) − 17)
Now, we can solve this differential equation using the method of variation of parameters.
Using formula (5) in Section 4.2,
e^(-P(x)) dx V₂ = V₁(x)
we can write the general solution as:
y₂ = c₁y₁ + c₂y₁ ∫ e^(-∫P(x)dx) dx
We can integrate e^(-∫P(x)dx) as follows:
∫ e^(-∫P(x)dx) dx = e^(-∫P(x)dx)
We need to find -∫P(x)dx. We have:
p(x) = −(1/x) − 4 sin(4 ln(x))/cos(4 ln(x))
So, -P(x) = ∫p(x)dx = −ln(x) + 4 ln(cos(4 ln(x)))
Therefore, e^(-∫P(x)dx) = x e^(-4 ln(cos(4 ln(x)))) = x cos^4( ln(x))
Now, we can write the second solution as:
y₂ = c₁x y cos(4 ln(x)) + c₂x y sin(4 ln(x))
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Consider the following differential equation. x′′+xx′−4x+x^3=0. By introducing a new variable y=x′, we set up a system of differential equations and investigate the behavior of its solution around its critical points (a,b). Which point is a unstable spiral point in the phase plane? A. (0,0) B. (1,3) C. (2,0) D. (−2,0)
To determine which point is an unstable spiral point in the phase plane for the given differential equation, we need to investigate the behavior of the solution around its critical points.
First, let's find the critical points by setting x' = 0 and x'' = 0 in the given differential equation. We are given the differential equation x'' + xx' - 4x + x^3 = 0.
Setting x' = 0, we get:
0 + x(0) - 4x + x^3 = 0
Simplifying the equation, we have:
x(0) - 4x + x^3 = 0
Next, setting x'' = 0, we get:
0 + x(0)x' - 4 + 3x^2(x')^2 + x^3x' = 0
Since we have introduced a new variable y = x', we can rewrite the equation as a system of differential equations:
x' = y
y' = -xy + 4x - x^3
Now, let's analyze the behavior of the solutions around the critical points (a, b). To do this, we need to find the Jacobian matrix of the system:
J = |0 1|
|-y 4-3x^2|
Now, let's evaluate the Jacobian matrix at each critical point:
For point (0,0):
J(0,0) = |0 1|
|0 4|
The eigenvalues of J(0,0) are both positive, indicating an unstable node.
Fopointsnt (1,3):
J(1,3) = |0 1|
|-3 1|
The eigenvalues of J(1,3) are both complex with a positive real part, indicating an unstable spiral point.
For point (2,0):
J(2,0) = |0 1|
|0 -eigenvalueslues lueslues of J(2,0) are both negative, indicating a stable node.
For point (-2,0):
J(-2,0) = |0 1|
|0 4|
The eigenvalues of J(-2,0) are both positive, indicatinunstablethereforebefore th hereherefthate point (1,3) is an unstable spiral point in the phase plane.
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Suppose that X and Y are independent random variables. If we know that E(X)=−5 and E(Y)=−2, determine the value of E(XY−6X). A. 40 B. 22 C. 10 D. −20 E. −2
The value of E(XY−6X) is 40.
To find the value of E(XY−6X), we can use the linearity of expectations. Since X and Y are independent random variables, the expected value of their product is equal to the product of their expected values.
E(XY) = E(X) * E(Y)
Given that E(X) = -5 and E(Y) = -2, we can substitute these values into the equation:
E(XY) = (-5) * (-2) = 10
Next, we need to calculate the expected value of -6X. Again, using the linearity of expectations:
E(-6X) = -6 * E(X)
Substituting the value of E(X) = -5:
E(-6X) = -6 * (-5) = 30
Now, we can find the expected value of the expression XY−6X by subtracting E(-6X) from E(XY):
E(XY−6X) = E(XY) - E(-6X) = 10 - 30 = -20
Therefore, the value of E(XY−6X) is -20.
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What are the differences between average and
instantaneous rates of change? Define
secant and tangent lines, and
explain how they are involved.
The average rate of change is the ratio of change in y-values to the change in x-values over a specific interval of time. The instantaneous rate of change is the rate of change at an exact point in time or space.
In calculus, secant lines are used to approximate a curve on a graph by drawing a line that intersects two points on the curve. On the other hand, a tangent line is a straight line that only touches a curve at one point and does not intersect it.
The average rate of change is used to estimate how quickly a function changes over a certain interval of time. In contrast, the instantaneous rate of change calculates the change at an exact moment or point. When we take the average rate of change over smaller and smaller intervals, the resulting values get closer to the instantaneous rate of change.
This is where the concept of tangent lines comes in. We use tangent lines to find the instantaneous rate of change of a function at a specific point. A tangent line touches a curve at a single point and represents the instantaneous rate of change at that point. On the other hand, a secant line is a line that intersects two points on a curve. It is used to approximate the curve of the function between the two points.
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