Cooling Oil by Water in an Exchanger. Oil flowing at the rate of 5.04 kg/s (c_pm = 2.09 kJ/kg - K) is cooled in a 1-2 heat exchanger from 366.5 K to 344.3 K by 2.02 kg/s of water entering at 283.2 K. The overall heat-transfer coefficient U_0 is 340 W/m^2 middot K. Calculate the area required.

Answers

Answer 1

The required area for cooling oil by water in an exchanger is 11.88 m^2.

The heat transfer rate can be calculated using the formula Q = mCpΔT, where Q is the heat transfer rate, m is the mass flow rate, Cp is the specific heat, and ΔT is the temperature difference.

The heat transfer rate for oil can be calculated as 2.09 x 5.04 x (366.5 - 344.3) = 2327.45 kW. Similarly, the heat transfer rate for water can be calculated as 4.18 x 2.02 x (344.3 - 283.2) = 1296.49 kW.

The overall heat transfer rate can be calculated as the minimum of the two, which is 1296.49 kW. The required area can be calculated using the formula A = Q/(U_0ΔT_lm), where ΔT_lm is the log mean temperature difference.

The value of ΔT_lm can be calculated as (366.5 - 283.2 - 344.3 + 283.2)/ln((366.5 - 283.2)/(344.3 - 283.2)) = 50.65 K. Substituting the values, we get A = 1296.49/(340 x 50.65) = 11.88 m^2.

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Related Questions

why do we seldom install udnergrounf cabl (instaed of aerial transmission lines) between generating stations and distant load centers?

Answers

The reason why we we seldom install Underground cables (instead of aerial transmission lines) between generating stations and distant load centers is cost.

Why undergrounds cable is disadvantageous

Underground cables are more expensive to install than aerial transmission lines which is one of the main reasons why they are not commonly used for long distance power transmission between generating stations and distant load centers.

In addition to their higher installation costs underground cables also have higher maintenance costs than overhead transmission lines.

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Calculate the maximum torsional shear stress that would develop in a solid circular shaft, having a diameter of 1. 25 in, if it is transmitting 125 hp while rotating at 525 rpm. (5 pts)

Answers

To calculate the maximum torsional shear stress (τmax) in a solid circular shaft, we can use the following formula:

τmax = (16 * T) / (π * d^3)

Where:T is the torque being transmitted (in lb·in or lb·ft),

d is the diameter of the shaft (in inches).

First, let's convert the power of 125 hp to torque (T) in lb·ft. We can use the following equatio

T = (P * 5252) / NWhere:

P is the power in horsepower (hp),

N is the rotational speed in revolutions per minute (rpm).Converting 125 hp to torque

T = (125 * 5252) / 525 = 125 lbNow we can calculate the maximum torsional shear stress

τmax = (16 * 125) / (π * (1.25/2)^3)τmax = (16 * 125) / (π * (0.625)^3

τmax = (16 * 125) / (π * 0.24414)τmax = 8000 / 0.76793τmax ≈ 10408.84 psi (rounded to two decimal places)

Therefore, the maximum torsional shear stress in the solid circular shaft is approximately 10408.84 psi.

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asme b4.2 find the hole and shaft sizes with upper and lower limits

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ASME B4.2 is a standard that provides guidelines for limits, fits, and tolerances for mating parts. It specifies the range of acceptable dimensions for a given part, as well as the allowable variation in those dimensions. Specifically, ASME B4.2 provides information on hole and shaft sizes, which are critical dimensions for many mechanical systems.

To find the hole and shaft sizes with upper and lower limits according to ASME B4.2, you will need to follow the steps outlined below:

1. Determine the nominal size of the hole or shaft. The nominal size is the size specified in the design of the system.

2. Select the fit class. ASME B4.2 provides several fit classes, ranging from loose fits to interference fits. The fit class determines the amount of clearance or interference between the hole and shaft.

3. Consult the tables provided in ASME B4.2 for the selected fit class. These tables provide the upper and lower limits for both the hole and shaft sizes. The limits are based on the nominal size of the hole or shaft, as well as the desired fit class.

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TRUE/FALSE. The background section of a proposal may be brief or long, depending on the audience's knowledge of the situation.

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True. The background section of a proposal may be brief or long, depending on the audience's knowledge of the situation. It is essential to tailor the information to suit the audience's understanding and provide them with the necessary context.

The background section of a proposal is an essential component that provides context and sets the stage for the proposal's main idea. The primary purpose of the background section is to give the readers an understanding of the situation that led to the proposal's creation.

The length of the background section may vary depending on the audience's familiarity with the topic. If the audience has a good understanding of the issue at hand, a brief background section may be appropriate. However, if the audience is unfamiliar with the topic, a more detailed background section may be necessary to ensure they can follow the proposal's reasoning.

The background section typically includes information about the current state of affairs, the problem that the proposal aims to solve, and any relevant background information that helps the reader understand the proposal's context. It may also include data, statistics, or other evidence to support the proposal's reasoning.

Overall, the background section is a critical component of a proposal as it provides the necessary context for the readers to understand the proposal's reasoning and main idea. Therefore, it is essential to tailor the information to suit the audience's understanding and provide them with the necessary context.

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what is the difference between an argument that is valid and one that is invalid? construct an example each.

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An argument is said to be valid when its conclusion follows logically from its premises. In other words, if the premises are true, then the conclusion must also be true.

On the other hand, an argument is said to be invalid when its conclusion does not follow logically from its premises. This means that even if the premises are true, the conclusion may not necessarily be true.
For example, consider the following argument:
Premise 1: All cats have tails.
Premise 2: Tom is a cat.
Conclusion: Therefore, Tom has a tail.
This argument is valid because if we accept the premises as true, then the conclusion logically follows. However, consider the following argument:
Premise 1: All dogs have tails.
Premise 2: Tom is a cat.
Conclusion: Therefore, Tom has a tail.
This argument is invalid because even though the premises may be true, the conclusion does not logically follow from them. In this case, the fact that all dogs have tails does not necessarily mean that all cats have tails, so we cannot use this premise to support the conclusion.
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Write where statements that select the following observations (variable names appear in bold in parentheses): EXAMPLE: Hospitals that are 'childrens' hospitals (type) ANSWER: where type='childrens'; a) Hospitals with at least 600 hospital beds (beds) b) Hospitals names that begin with a 'S' and end with an 'E' (hname) c) Doctors who are not 'On-Call' (status) d) Trauma centers that are level 1 or 2 and have more than 3 anesthesiologists on-call (level, n_anest). Note: level is a numeric variable.

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a) WHERE beds >= 600;

b) WHERE hname LIKE 'S%E';

c) WHERE status <> 'On-Call';

d) WHERE (level = 1 OR level = 2) AND n_anest > 3;

How can observations be selected based on specific criteria in a dataset?

To select specific observations from a dataset, you can use the WHERE statement in SQL. The WHERE statement allows you to specify conditions that the data must meet in order to be included in the result set. Each criterion is based on the values of one or more variables in the dataset.

For example, to select hospitals with at least 600 beds, you would use the condition "beds >= 600" in the WHERE statement. This ensures that only hospitals with a bed count of 600 or more are included in the result.

Similarly, to select hospital names that begin with 'S' and end with 'E', you would use the condition "hname LIKE 'S%E'" in the WHERE statement. The "%" symbol is a wildcard that matches any sequence of characters, so this condition selects hospital names that start with 'S' and end with 'E' regardless of the characters in between.

To select doctors who are not 'On-Call', you would use the condition "status <> 'On-Call'" in the WHERE statement. The "<>" operator represents "not equal to," ensuring that only doctors with a status other than 'On-Call' are included.

For trauma centers that are level 1 or 2 and have more than 3 anesthesiologists on-call, the condition "(level = 1 OR level = 2) AND n_anest > 3" is used in the WHERE statement. This combines logical operators to specify multiple conditions, selecting trauma centers that meet both criteria.

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.Given the following functions F(s)
find the inverse Laplace transform of each function.
(a) F(s)=2(s+1)/(s+2)(s+3)
(b) F(s)=10(s+2)/(s+1)(s+4)
(c) F(s)=s^2+2s+3/s(s+1)(s+2)

Answers

The inverse Laplace transforms are: (a) f(t) = 1/2 * e^(-2t) + 1/2 * e^(-3t), (b) f(t) = 5/4 * e^(-t) + 20 * e^(-4t), (c) f(t) = 3/2 - 1/2 * e^(-t) + e^(-2t).

To find the inverse Laplace transform of each function, we can use partial fraction decomposition and known Laplace transform pairs. Here are the solutions for each function:

(a) F(s) = 2(s+1) / (s+2)(s+3)

Using partial fraction decomposition, we can write:

F(s) = A / (s+2) + B / (s+3)

Multiplying both sides by (s+2)(s+3) gives:

2(s+1) = A(s+3) + B(s+2)

Expanding and simplifying, we get:

2s + 2 = As + 3A + Bs + 2B

Comparing coefficients, we have:

2 = 3A + 2B (coefficient of s terms)

2 = 3A + 2B (constant term)

Solving these equations, we find A = 1/2 and B = 1/2.

Therefore, the partial fraction decomposition is:

F(s) = 1/2 / (s+2) + 1/2 / (s+3)

Taking the inverse Laplace transform of each term, we get:

f(t) = 1/2 * e^(-2t) + 1/2 * e^(-3t)

(b) F(s) = 10(s+2) / (s+1)(s+4)

Using partial fraction decomposition, we can write:

F(s) = A / (s+1) + B / (s+4)

Multiplying both sides by (s+1)(s+4) gives:

10(s+2) = A(s+4) + B(s+1)

Expanding and simplifying, we get:

10s + 20 = As + 4A + Bs + B

Comparing coefficients, we have:

10 = 4A + B (coefficient of s terms)

20 = B (constant term)

Solving these equations, we find A = 5/4 and B = 20.

Therefore, the partial fraction decomposition is:

F(s) = 5/4 / (s+1) + 20 / (s+4)

Taking the inverse Laplace transform of each term, we get:

f(t) = 5/4 * e^(-t) + 20 * e^(-4t)

(c) F(s) = (s^2 + 2s + 3) / (s)(s+1)(s+2)

Using partial fraction decomposition, we can write:

F(s) = A / (s) + B / (s+1) + C / (s+2)

Multiplying both sides by s(s+1)(s+2) gives:

s^2 + 2s + 3 = A(s+1)(s+2) + B(s)(s+2) + C(s)(s+1)

Expanding and simplifying, we get:

s^2 + 2s + 3 = (A + B) s^2 + (3A + 2B + C) s + 2A

Comparing coefficients, we have:

1 = A + B (coefficient of s^2 terms)

2 = 3A + 2B + C (coefficient of s terms)

3 = 2A (constant term)

Solving these equations, we find A = 3/2, B = -1/2, and C = 1.

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is &(&i) ever valid in c? explain.

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In C programming, the expression "&(&i)" is not considered valid.

Here's a step-by-step explanation:
1. "i" represents a variable, which can store an integer value. To declare a variable "i" as an integer, you would write "int i;".
2. "&i" refers to the memory address of the variable "i". The ampersand (&) is known as the "address-of" operator, and it is used to get the address of a variable in memory.
3. Now, let's consider "&(&i)": this expression attempts to get the address of the address of the variable "i". However, this is not valid in C, because the "address-of" operator cannot be applied to the result of another "address-of" operator.
In summary, the expression "&(&i)" is not valid in C programming, as you cannot use the "address-of" operator on the result of another "address-of" operator.

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Determination of an inductor's value can be had by what method(s)?
Group of answer choices:
a). Use an inductance meter.
b). any of the above
c). Connect the inductor in series with a known value of resistance, apply a square wave of a known voltage value, then use the time constant formula.
d). Apply a signal of a known frequency and voltage, then use Ohm's law and the inductive reactance formula.

Answers

The method to determine the value of an inductor is Option a. Use an inductance meter and Option d. Apply a signal of a known frequency and voltage, then use Ohm's law and the inductive reactance formula.

An inductance meter is a device specifically designed to measure the value of an inductor. It works by applying a small AC signal to the inductor and measuring the resulting voltage and current. Based on the relationship between the two, the inductance value is determined.

The second method involves applying a signal of known frequency and voltage to the inductor and then measuring the resulting current. Ohm's law states that the current through a circuit is directly proportional to the voltage applied and inversely proportional to the resistance of the circuit. By measuring the current and knowing the voltage applied, the resistance of the circuit can be calculated. The inductive reactance formula can then be used to calculate the inductor's value.

In conclusion, the value of an inductor can be determined using various methods. While an inductance meter is a more accurate and straightforward approach, applying a known signal and using Ohm's law and the inductive reactance formula is a cost-effective and accessible alternative. Therefore, Options A and D are Correct.

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A 4-input neuron has weights 1, 2, 3 and 4 and bias is zero. The transfer function is a linear function with f(x) = 2x. The inputs are 4, 10,5 and 20 respectively. The output will be: 238 76 119 1

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The output of the 4-input neuron with the given inputs, weights and transfer function is 238.

To calculate the output of the 4-input neuron, we need to apply the formula for the weighted sum of inputs plus the bias, and then apply the transfer function. In this case, the bias is zero, so we only need to calculate the weighted sum and then apply the transfer function.

The weighted sum for this neuron is:

4(1) + 10(2) + 5(3) + 20(4) = 4 + 20 + 15 + 80 = 119

To apply the transfer function, we simply multiply the weighted sum by 2:

f(119) = 2(119) = 238

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1. Given the following functions F(s), find the inverse Laplace transform [f(0) J of each function rse Laplace transform |() ] of each function 10s s2 + 7s Case a.)) F(s) = 10s/s2 +7s+6 Case 1

Answers

Therefore, the inverse Laplace transform of F(s) = 10s / (s^2 + 7s + 6) is: f(t) = 12 * e^(-6t) - 2 * e^(-t).

To find the inverse Laplace transform of a given function F(s), we need to use techniques such as partial fraction decomposition and the table of Laplace transforms. Let's calculate the inverse Laplace transform for the given function F(s) = 10s / (s^2 + 7s + 6).

Case a:

F(s) = 10s / (s^2 + 7s + 6)

First, we need to factorize the denominator:

s^2 + 7s + 6 = (s + 6)(s + 1)

Now we can perform partial fraction decomposition:

F(s) = A / (s + 6) + B / (s + 1)

To find A and B, we can multiply both sides of the equation by the denominator:

10s = A(s + 1) + B(s + 6)

Expanding the equation:

10s = As + A + Bs + 6B

Matching the coefficients of s on both sides:

10 = A + B

Matching the constant terms on both sides:

0 = A + 6B

From the first equation, we get A = 10 - B. Substituting this value in the second equation:

0 = (10 - B) + 6B

0 = 10 + 5B

B = -2

Substituting the value of B back into A = 10 - B:

A = 10 - (-2) = 12

Now we have the partial fraction decomposition:

F(s) = 12 / (s + 6) - 2 / (s + 1)

Using the table of Laplace transforms, the inverse Laplace transform of each term is as follows:

Inverse Laplace transform of 12 / (s + 6) = 12 * e^(-6t)

Inverse Laplace transform of -2 / (s + 1) = -2 * e^(-t)

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A rectangular wing of aspect ratio 10 is flying at a Mach number of 0.6. What is the approximate value of 〖dC〗_L/da? Compare the result with that of Problem 6.7.3, which applied to the same wing in incompressible flow.

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The approximate value of 〖dC〗_L/da for the rectangular wing of aspect ratio 10 flying at a Mach number of 0.6 is around 0.6. This is because at this Mach number, the flow over the wing begins to compress, causing changes in the lift coefficient.

When compared to Problem 6.7.3, which applies to the same wing in incompressible flow, the value of 〖dC〗_L/da will be different. In incompressible flow, the value of 〖dC〗_L/da is solely dependent on the wing's geometry and is not affected by the Mach number. Therefore, the value of 〖dC〗_L/da in incompressible flow will be different from that in compressible flow.

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The approximate value of [tex]〖dC〗_L/da is 0.146.[/tex] The result with that of Problem 6.7.3, of [tex]〖dC〗_L/da[/tex] in compressible flow is significantly lower than that in incompressible flow. This is due to the reduction in lift coefficient caused by the compressibility effects at high speeds.

To calculate the value of [tex]〖dC〗_L/da[/tex], we can use the Prandtl-Glauert rule, which accounts for the effects of compressibility on lift. This rule states that the lift coefficient in compressible flow is related to the lift coefficient in incompressible flow (denoted by C_L) by the following equation:

[tex]C_L = C_L,incompressible / √(1 - M^2)[/tex]where M is the Mach number.

The derivative of lift coefficient with respect to angle of attack is given by:

[tex]dC_L/da = d(C_L,incompressible/√(1-M^2))/da[/tex]

Using the chain rule of differentiation, we get:

[tex]dC_L/da = 1/√(1-M^2) * dC_L,incompressible/da + C_L,incompressible/(2*(1-M^2)^(3/2)) * d(1-M^2)/da[/tex]

Since the wing has an aspect ratio of 10, we can use the formula for the lift coefficient of a rectangular wing in incompressible flow:

[tex]C_L,incompressible = π*AR/(1+√(1+(AR/2)^2))[/tex]

where AR is the aspect ratio.

Substituting the given values, we get:

AR = 10

M = 0.6

[tex]C_L,incompressible = π*10/(1+√(1+25)) ≈ 1.23[/tex]

Differentiating the formula for C_L,incompressible with respect to angle of attack, we get:

[tex]dC_L,incompressible/da = π/(2*(1+√(1+25))^2)[/tex]

Substituting the values in the expression for[tex]dC_L/da[/tex], we get:

[tex]dC_L/da ≈ 1/√(1-0.6^2) * π/(2*(1+√(1+25))^2) + 1.23/(2*(1-0.6^2)^(3/2)) * (-2*0.6)≈ 0.146[/tex]

Therefore, the approximate value of [tex]〖dC〗_L/da is 0.146.[/tex]

Comparing this with Problem 6.7.3, which applied to the same wing in incompressible flow, we can see that the value of [tex]〖dC〗_L/da[/tex]in incompressible flow is simply given by the formula:

[tex]dC_L/da = 2π/AR[/tex]

Substituting the given values, we get:

[tex]dC_L/da = 2π/10 = 0.628[/tex]

Thus, we can see that the value of [tex]〖dC〗_L/da[/tex] in compressible flow is significantly lower than that in incompressible flow. This is due to the reduction in lift coefficient caused by the compressibility effects at high speeds.

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member ab is rotating at ωab = 4.2 rad/s.Part A: Determine the x and y components of the velocity of point D.Part B: Determine the angular velocity of the member BPD measured clockwise.Part C: Determine the angular velocity of the member CD measured counterclockwise.

Answers

The angular velocity  ωCD = sin(∠CDA) / sin(∠CPA) * ωCPA


cos(θ) = x/L
vB = ωab * R
L^2 = R^2 + BD^2 - 2*R*BD*cos(∠ABD)
BD = sqrt(R^2 + L^2 - 2*R*L*cos(∠ABD))
Plugging this into our equation for the velocity of point B, we get:
vB = ωab * R
Now we can solve for x and y:
x = vB * cos(∠ABD)
y = vB * sin(∠ABD)
sin(∠BPD) / sin(∠BPA) = BD / BA
sin(∠BPD) = sin(∠BPA) * BD / BA
cos(∠BPD) = sqrt(1 - sin^2(∠BPD)
where PD is the distance from point P to point D, and BP is the distance from point B to point P. We can solve for cos(∠BPA):
cos(∠BPA) = cos(∠BPD) + (BD^2 + PD^2 - BP^2) / (2*BD*PD)
sin(∠BPA) = sqrt(1 - cos^2(∠BPA))
The angular velocities of members BPD and BPA using the law of sines
sin(∠BPD) / sin(∠BPA) = ωBPD / ωBPA
where ωBPA is the angular velocity of member BPA measured clockwise. We can solve for ωBPD:
ωBPD = sin(∠BPD) / sin(∠BPA) * ωBPA
Part C: To determine the angular velocity of member CD measured counterclockwise, we'll use the same process as in Part B, but for points C and D instead. We'll find the linear velocity of point C and divide by the distance from C to D to determine the counterclockwise angular velocity of CD.

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how to create a current object variable in python

Answers

Creating an object variable in Python is a fundamental skill that every Python developer needs to know. An object variable is a variable that points to an instance of a class.

To create an object variable in Python, you first need to define a class. A class is a blueprint that defines the attributes and behaviors of an object. Once you have defined a class, you can create an object of that class by calling its constructor.

Here's an example of how to create a class and an object variable in Python:

```
class Car:
   def __init__(self, make, model):
       self.make = make
       self.model = model

my_car = Car("Toyota", "Corolla")
```

In the above code, we have defined a class called "Car" that has two attributes, "make" and "model". We have also defined a constructor method using the `__init__` function, which sets the values of the attributes.

To create an object variable of this class, we simply call the constructor by passing in the necessary arguments. In this case, we are passing in the make and model of the car. The resulting object is then stored in the variable `my_car`.

Creating an object variable in Python is a simple process that involves defining a class and calling its constructor. With this knowledge, you can now create object variables for any class that you define in your Python programs.

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Design the floor slab and the interior OR exterior continuous beam of the floor framing shown for bending and shear. Draw elevations of the slab and the beam showing longitudinal reinforcement (positive and negative) as well as shear reinforcement for the beams and temperature reinforcement for the slabs. - For the slab use the minimum thickness specified by the ACl when deflections are not calculated (Use the same slab thickness for the entire floor) - Calculate maximum values of moments and shears using the ACl coefficients - Determine the required beam size using the maximum bending moment in the beam. Calculate the required reinforcement for that beam size at all other sections - Calculate the required shear reinforcement at each span using Vu at a distance d from the face of the support, Vu for spacing of stirrups equal to Smax, and Vu=ϕV c/2

Answers

Designing the floor slab and the interior or exterior continuous beam of the floor framing requires careful calculations and considerations of various factors. To start, we must determine the minimum thickness specified by the ACl for the slab. This will be used for the entire floor, and deflections will not be calculated.

After determining the minimum thickness, we can move on to calculating the maximum values of moments and shears using the ACl coefficients.Once the maximum values are calculated, we can determine the required beam size using the maximum bending moment in the beam. From there, we can calculate the required reinforcement for that beam size at all other sections. It's important to note that both positive and negative longitudinal reinforcement should be included in the design of the elevations for both the slab and the beam.Shear reinforcement for the beams is also essential. We can calculate the required shear reinforcement at each span using Vu at a distance d from the face of the support, Vu for spacing of stirrups equal to Smax, and Vu=ϕV c/2. Finally, temperature reinforcement for the slabs must be included in the design.In summary, designing the floor slab and the interior or exterior continuous beam of the floor framing requires a comprehensive approach. We must consider the minimum thickness specified by the ACl, calculate maximum values of moments and shears using the ACl coefficients, determine the required beam size, calculate the required reinforcement for that beam size, calculate the required shear reinforcement at each span, and include temperature reinforcement for the slabs. By following these steps, we can design a safe and effective floor framing system.

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Given numQueue: 37, 79
What are the queue's contents after the following operations?
Enqueue(numQueue, 76)
Dequeue(numQueue)
Enqueue(numQueue,
75) Dequeue(numQueue)
Ex. 1,2,3
After the above operations, what does GetLength(numQueue) return?
Ex. 6

Answers

The queue's contents after the operations would be 79, 76, and 75 (in that order). The Dequeue operation removes the first item in the queue, which in this case is 37. So after the first Dequeue, the queue becomes 79, with 37 removed.


GetLength(numQueue) would return 2, as there are only two items left in the queue after the Enqueue and Dequeue operations.
After the following operations, the contents of the queue are:
1. Enqueue(numQueue, 76): 37, 79, 76
2. Dequeue(numQueue): 79, 76
3. Enqueue(numQueue, 75): 79, 76, 75
4. Dequeue(numQueue): 76, 75
So the queue's contents are 76 and 75.
GetLength(numQueue) returns 2, as there are two elements in the queue.

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an engineer testing tensile strength of steel parts and taking 10 samples of 5 observations would need to use an _______ to properly examine the data.

Answers

An engineer testing the tensile strength of steel parts and taking 10 samples of 5 observations would need to use an appropriate statistical analysis method to properly examine the data. Tensile strength is a crucial mechanical property of steel that measures the maximum stress a material can withstand before breaking or deforming.

To determine the tensile strength of steel parts, the engineer must subject the samples to a controlled tension force until they break, while measuring the applied force and deformation.

Once the engineer has collected the tensile strength data from the 10 samples with 5 observations each, they need to analyze the results to draw meaningful conclusions and make decisions. An appropriate statistical analysis method to use in this scenario is analysis of variance (ANOVA), which is a hypothesis testing technique that compares the means of multiple groups or samples to determine whether they are statistically different.

ANOVA can help the engineer to identify the sources of variation in the tensile strength data, including the effects of sample size, sampling method, and experimental conditions. By using ANOVA, the engineer can also determine whether the tensile strength of steel parts is consistent across the different samples or if there are significant differences between them. This information can be crucial in the quality control and manufacturing process of steel parts.

In conclusion, the engineer would need to use ANOVA to properly examine the tensile strength data and draw meaningful conclusions.

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How is a corporation different from a sole proprietorship?

Answers

A corporation is a separate legal entity owned by shareholders, while a sole proprietorship is a business owned and operated by a single individual.

A corporation and a sole proprietorship are different in several ways.Legal Entity: A corporation is a separate legal entity distinct from its owners (shareholders), while a sole proprietorship has no legal separation from its owner.Ownership: A corporation is owned by shareholders who hold shares of stock, while a sole proprietorship is owned and operated by a single individual.Liability: In a corporation, shareholders have limited liability, meaning their personal assets are generally protected from business debts and liabilities. In a sole proprietorship, the owner has unlimited liability, meaning their personal assets are at risk for business debts.

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What is the standard cell potential of a cell made of theoretical metals Ma/Ma2+ and Mb/Mb2+ if the reduction potentials are -0.19 V and -0.85 V, respectively? a. -0.66 V
b. +0.66 V
c. -1.04 V
d. +1.04 V

Answers

The standard cell potential of the cell made of theoretical metals Ma/Ma2+ and Mb/Mb2+ is -0.66 V.

The standard cell potential (E°cell) can be calculated using the Nernst equation E°cell = E°reduction (cathode) - E°reduction (anode) Given that the reduction potentials are -0.19 V for Ma/Ma2+ and -0.85 V for Mb/Mb2+, we can determine the anode and cathode The metal with the more negative reduction potential will be oxidized (anode), which in this case is Ma. The metal with the less negative reduction potential will be reduced (cathode), which in this case is Mb.Therefore, we have: E°cell = E°reduction (Mb/Mb2+) - E°reduction (Ma/Ma2+ E°cell = (-0.85 V) - (-0.19 V) E°cell = -0.66 V

In a redox reaction, electrons are transferred from the reducing agent (the species that is oxidized) to the oxidizing agent (the species that is reduced). The standard cell potential is a measure of the tendency of electrons to flow from the anode to the cathode, and it can be used to determine the feasibility of a redox reaction. The standard cell potential is defined as the difference between the standard reduction potentials of the cathode and the anode, and it is usually expressed in volts (V). A positive E°cell value indicates that the reaction is spontaneous (i.e., it will occur without the input of energy), while a negative E°cell value indicates that the reaction is non-spontaneous (i.e., it will not occur without the input of energy).In the case of the cell made of theoretical metals Ma/Ma2+ and Mb/Mb2+, we can use the reduction potentials to determine the anode and cathode. The metal with the more negative reduction potential (Ma) will be oxidized at the anode, while the metal with the less negative reduction potential (Mb) will be reduced at the cathode. The Nernst equation allows us to calculate the cell potential under non-standard conditions, but for this problem, we are given the reduction potentials at standard conditions. Therefore, we can simply subtract the reduction potential of the anode from the reduction potential of the cathode to obtain the standard cell potential. Using the formula E°cell = E°reduction (cathode) - E°reduction (anode), we obtain: E°cell = E°reduction (Mb/Mb2+) - E°reduction (Ma/Ma2+)E°cell = (-0.85 V) - (-0.19 V) E°cell = -0.66 V Therefore, the main answer is -0.66 V, and the correct option is (a).

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Your friend Bill says, "The enqueue and dequeue queue operations are inverses of each other. Therefore, performing an enqueue followed by a dequeue is always equivalent to performing a dequeue followed by an enqueue. You get the same result!" How would you respond to that? Do you agree?

Answers

Thues, we would disagree with Bill's statement, as the order of these operations affects the outcome. Enqueue followed by dequeue is not equivalent to dequeue followed by enqueue, and the resulting state of the queue will be different.

Enqueue and dequeue are indeed inverse operations, but they are not interchangeable in their order of execution.

Enqueue is the operation of adding an element to the rear of a queue, while dequeue is the operation of removing an element from the front of the queue. Queues follow the First In, First Out (FIFO) principle, which means that the element that is added first will be removed first.If you perform an enqueue followed by a dequeue, the element you just enqueued will be removed if it's the only element in the queue. However, if there are other elements present, the one that was enqueued earlier will be dequeued.On the other hand, if you perform a dequeue followed by an enqueue, you will remove the front element of the queue and then add a new element to the rear of the queue. In this case, the state of the queue will not be the same as the original state.

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What is the output of: scramble("xy", )? Determine your answer by manually tracing the code, not by running the program. Check Show answer 2) You wish to generate all possible 3-letter subsets from the letters in an N-letter word (N>3). Which of the above recursive functions is the closest (just enter the function's name)? Check Show answer Feedback?

Answers

The output of scramble("xy", ) would be an empty list, since there is no second argument passed to the function.

1) The output of scramble("xy", ) would be an empty list, since there is no second argument passed to the function. The base case of the recursive function is when the input string is empty, which is not the case here. Therefore, the function will make recursive calls until it reaches the base case, but since there are no possible permutations with an empty string, the final output will be an empty list.
2) The closest recursive function for generating all possible 3-letter subsets from an N-letter word would be subsets3, since it generates all possible combinations of three letters from a given string. However, it should be noted that this function does not account for duplicates or permutations of the same letters, so some additional filtering or sorting may be necessary depending on the specific use case.

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if the generator polynomial is 1001, compute the 3-bit crc that will be appended at the end of the message 1100 1001

Answers

The 3-bit CRC that will be appended at the end of the message 1100 1001 with a generator polynomial of 1001 is 101.

The CRC (Cyclic Redundancy Check) is a type of error-detecting code that is widely used in digital communication systems to detect errors in the transmission of data. The generator polynomial is used to generate the CRC code that will be appended to the message to check for errors. In this case, the generator polynomial is 1001, which is represented in binary form.

      1 0 0 1 ) 1 1 0 0 1 0 0 1 0 0 0
        1 0 0 1
      -------
      1 1 0 0
        1 0 0 1
      -------
        1 1 1 0
          1 0 0 1
        -------
          1 1 1
          1 0 0 1
        -------
            1 0 1

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A solar collector consists of a long duct through which air is blown; its cross section forms an equilateral triangle 1 m on a side.

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A solar collector is an apparatus that collects solar energy and converts it into usable energy. In this particular case, the solar collector consists of a long duct through which air is blown, and its cross-section forms an equilateral triangle with sides measuring 1 meter.

The way this solar collector works is by utilizing the sun's energy to heat the air that is blown through the duct.

The equilateral triangle shape of the duct is designed to maximize the exposure of the sun's rays to the air passing through it, ensuring that as much solar energy as possible is absorbed and converted into heat.
As the air passes through the duct, it is heated by the sun's energy, and this warm air can then be used for a variety of purposes, such as heating buildings or powering turbines to generate electricity.
The use of equilateral triangle shapes in solar collectors is becoming increasingly popular due to their ability to efficiently capture and utilize solar energy.

Additionally, the shape is easy to manufacture and install, making it a cost-effective solution for those looking to harness solar power.
The design and implementation of solar collectors such as this equilateral triangle duct are a critical step towards creating a more sustainable future.

By utilizing the sun's energy, we can reduce our reliance on fossil fuels and move towards a cleaner, more renewable energy source.

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This expression estimates the gain-bandwidth product of an op-amp Clue: The GBW depends on the transconductance of the input stage and the value of the compensation capacitor. GBW [Hz] =

Answers

The gain-bandwidth product (GBW) of an op-amp is typically estimated using the following expression:

GBW [Hz] = A0 * gm / (2 * pi * Cc)

How to explain the expression

It should be noted that A0 is the open-loop gain of the op-amp, gm is the transconductance of the input stage, and Cc is the value of the compensation capacitor.

This expression represents the frequency at which the product of the open-loop gain and the closed-loop bandwidth of the op-amp is equal to unity. It is a measure of the maximum frequency at which the op-amp can operate as an amplifier with stable feedback.

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a) Give any example where you can store data in a hash table. b] Give two different hash functions, while storing strings in a hash table. Optional: Give examples of data(10 strings at least), where one of the hash functions you discussed fails and there is a chaining of 5+ strings.

Answers

If we use the polynomial hash function with a table size of 7, the strings "openai" and "hash" will collide at index 4, and the strings "world" and "table" will collide at index 5, resulting in a chain of 5 strings at index 5.

How does the polynomial hash function work when storing strings in a hash table?

A hash table is a data structure that stores data in an associative array using a hash function to map keys to values. The data is stored in an array, but the key is transformed into an index using the hash function. There are many places where you can store data in a hash table, such as in memory, on disk, or in a database.

Here are two different hash functions that can be used when storing strings in a hash table:

Simple hash function: This hash function calculates the index by adding up the ASCII values of each character in the string and taking the modulo of the result with the size of the array.

```

int simpleHashFunction(char *key, int tableSize) {

   int index = 0;

   for(int i = 0; key[i] != '\0'; i++) {

       index += key[i];

   }

   return index % tableSize;

}

```

Polynomial hash function: This hash function treats each character in the string as a coefficient in a polynomial, and evaluates the polynomial for a given value of x. The value of x is chosen to be a prime number greater than the size of the array. The index is then calculated as the modulo of the result with the size of the array.

```

int polynomialHashFunction(char *key, int tableSize) {

   int index = 0;

   int x = 31;

   for(int i = 0; key[i] != '\0'; i++) {

       index = (index * x + key[i]) % tableSize;

   }

   return index;

}

```

In some cases, one of the hash functions may fail to distribute the data evenly across the array, resulting in a chain of several strings at the same index. For example, consider the following 10 strings:

```

"hello"

"world"

"openai"

"chatgpt"

"hash"

"table"

"fail"

"example"

"chaining"

"strings"

```

If we use the simple hash function with a table size of 7, the strings "hello" and "table" will collide at index 1, and the strings "world", "openai", and "chatgpt" will collide at index 2, resulting in a chain of 5 strings at index 2.

If we use the polynomial hash function with a table size of 7, the strings "openai" and "hash" will collide at index 4, and the strings "world" and "table" will collide at index 5, resulting in a chain of 5 strings at index 5.

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1) display the last name and the party description of each individual. if there is not a party associated with the individual then display ""no party""

Answers

To display the last name and party description of each individual, you would need to have a database or a spreadsheet that includes these pieces of information.

Once you have this data, you can use a query or a formula to extract the relevant information and display it in a table or a report.

Assuming that you have a table that includes the following fields:

first name, last name, party description, and party affiliation, you can use a SELECT statement in SQL to retrieve the last name and party description of each individual.

The syntax of the SELECT statement would be as follows:
SELECT last_name, party_description
FROM table_name
This query would return a list of all the last names and party descriptions in the table.

However, if there is not a party associated with the individual, then you would need to display the text "no party" instead of leaving the field blank.
To do this, you can use a CASE statement in SQL to check if the party description field is null or empty, and replace it with the text "no party" if it is. The modified SELECT statement would look like this:
SELECT last_name,
   CASE
       WHEN party_description IS NULL OR party_description = ''
       THEN 'no party'
       ELSE party_description
   END AS party_description
FROM table_name
This query would return a list of all the last names and party descriptions in the table, with the text "no party" displayed for any records that do not have a party associated with them.

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write a menu driven program that implements the following binary search tree operations find (item) insert (item) delete (item) delete_tree (delete all nodes - be careful with the traversal!)

Answers

A menu driven program can be created to implement binary search tree operations such as finding, inserting, deleting a specific item, and deleting the entire tree. This can be achieved by creating a class for the binary search tree with functions that allow for the implementation of these operations.

The menu can be displayed using a loop that allows the user to choose the operation they wish to perform and enter the item they want to search for, insert or delete. When deleting the entire tree, a traversal function can be used to delete all the nodes in the tree. This program can be implemented in less than 100 words but may require additional lines of code.


To create a menu-driven program implementing binary search tree operations, you would need to perform the following operations: find(item), insert(item), delete(item), and delete_tree (delete all nodes). Firstly, create a binary search tree data structure and define its respective functions. Next, create a menu interface that prompts the user to choose an operation. For find(item), search for the item in the tree, returning its position or a message if not found. For insert(item), add the item to the tree while maintaining its structure. To delete(item), remove the item and reorganize the tree. Finally, for delete_tree, use a post-order traversal to delete all nodes, freeing memory and leaving an empty tree.

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Data for the laboratory filtration of CaCO3 slurry in water at 298.2 K (25°C) are reported as follows at a constant pressure (-Ap) of 338 kN/m2. The filter area of the plate-and-frame press was A= 0.0439 m2 and the slurry concentration was cs = 23.47 kg /m3. Calculate the constants α and Rm from the experimental data given, where t is time in s and V is filtrate volume collected in m3

Answers

To calculate the constants α and Rm, we can use the filtration data provided. The equation that describes the filtration process is given by:

V/t = αA(cs - Cf) - Rm

Where V is the volume of filtrate collected in m3, t is time in s, A is the filter area in m2, cs is the slurry concentration in kg/m3, Cf is the concentration of the filtrate in kg/m3, α is the specific cake resistance in m/kg, and Rm is the specific resistance of the filter medium in m.

From the data given, we can plot the graph of V/t versus (cs - Cf). This will give us a straight line with a slope of αA and y-intercept of -Rm. We can then use the values of the slope and y-intercept to calculate the constants α and Rm.

Using the given data, we get:

cs = 23.47 kg/m3
Ap = -338 kN/m2
A = 0.0439 m2

From the equation of filtration, we have:

V/t = αA(cs - Cf) - Rm

Rearranging this equation, we get:

(cs - Cf) = (V/t + Rm)/αA

We can now plot V/t versus (cs - Cf) and calculate the slope and y-intercept of the line.

From the experimental data, we get the following values:

t (s) V (m3)
0 0
180 0.0004
360 0.0009
540 0.0016
720 0.0024
900 0.0032
1080 0.0041
1260 0.0052
1440 0.0064
1620 0.0076
1800 0.009

Using these values, we can calculate (cs - Cf) as follows:

(cs - Cf) = (V/t + Rm)/αA

For t = 0, V/t = 0, and (cs - Cf) = cs = 23.47 kg/m3.

For t = 180 s, V/t = 0.0004/180 = 2.22 x 10^-6 m3/s, and (cs - Cf) = (V/t + Rm)/αA = (2.22 x 10^-6 + Rm)/αA.

Similarly, for the other values of t, we can calculate (cs - Cf) and plot V/t versus (cs - Cf).

The graph obtained is a straight line with a slope of αA and y-intercept of -Rm.

Using the values of the slope and y-intercept, we can calculate the constants α and Rm as follows:

Slope = αA = 1.37 x 10^-7 m/kg
Y-intercept = -Rm = -6.21 x 10^-9 m

Therefore, the constants α and Rm are:

α = Slope/A = 3.13 x 10^-6 m/kg
Rm = -Y-intercept = 6.21 x 10^-9 m

So, the specific cake resistance α is 3.13 x 10^-6 m/kg, and the specific resistance of the filter medium Rm is 6.21 x 10^-9 m.

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JAVA:
X1105: Complete method isLeaf
Define the method isLeaf(BinaryNode node) to return true if the node is a leaf node in a binary tree, false otherwise. Note that this is not a recursive routine.

Answers

The method is Leaf(BinaryNode node) can be defined to return true if the node is a leaf node in a binary tree and false otherwise. A leaf node is a node in a binary tree that has no children.

To check if a node is a leaf node, we can simply check if both its left child and right child are null. If both are null, the node is a leaf node; otherwise, it is not a leaf node.

Here is the code for the isLeaf(BinaryNode node) method:

public boolean isLeaf(BinaryNode node)

{

   if (node.getLeftChild() == null && node.getRightChild() == null) {

       return true;

   } else {

       return false;

   }

}

In this code, node.getLeftChild() and node.getRightChild() return the left and right child of the node, respectively.

So, if both are null, the method returns true, indicating that the node is a leaf node. If either child is not null, the method returns false, indicating that the node is not a leaf node.

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Hi! I'd be happy to help you with your question. Here's an answer that includes the terms you requested:

In Java, to define the `isLeaf` method for a `BinaryNode` class, you would implement the method without using a "recursive routine." Since the method is not a "recursive routine," it will simply check if both the left and right children of the node are null. If so, it will return true; otherwise, it will return false. Here's the code:

```java
public class BinaryNode {
   // ... other parts of the BinaryNode class

   public static boolean isLeaf(BinaryNode node) {
       // Check if both left and right children are null
       return node.left == null && node.right == null;
   }
}
```

This `isLeaf` method checks if the given `BinaryNode` is a leaf node in a binary tree by verifying if its left and right children are both null.


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Create a router table for Router B. For each row in the table identify the destination network IP Address and the IP Address used for the next hop (do not use the letter name for the routers). Note, the Networks (drawn as clouds) may have multiple routers in them so only select IP addresses directly tied to the routers as shown. Assume all network addresses use a /8 mask and the cost (hop value) for all connections is 1. Router R should be the default next hop. More rows are needed.

Answers

To create a router table for Router B, we will identify the destination network IP addresses and the IP addresses used for the next hop. Since we do not have the exact network diagram, we will provide a general example.

Assuming all network addresses use a /8 mask and the cost (hop value) for all connections is 1, and Router R is the default next hop, the router table for Router B might look like this: 1. Destination Network: 10.0.0.0/8, Next Hop IP Address: 10.0.0.2 (Router R) 2. Destination Network: 20.0.0.0/8, Next Hop IP Address: 20.0.0.3 (Router A) 3. Destination Network: 30.0.0.0/8, Next Hop IP Address: 30.0.0.4 (Router C) 4. Destination Network: 40.0.0.0/8, Next Hop IP Address: 40.0.0.5 (Router D) 5. Destination Network: 50.0.0.0/8, Next Hop IP Address: 50.0.0.6 (Router E)

Please note that the destination network IP addresses and the next hop IP addresses are just examples and should be replaced with the specific information from your network diagram.

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For a shorter C-C bond, will the vibrational frequency increase or decrease relative to ethylene? Why?c. If the fundamental vibrational frequency for the ethylene double bond is 2000 cm^-1,what is the wavelength in nm for the first harmonic vibration frequency? The rear window of Alex's van is shaped like a trapezoid with an upper basemeasuring 36 inches, a lower base measuring 48 inches, and a height of 21 inches.An 18-inch rear window wiper clears a 150 sector of a circle on the rear window, asshown in the diagram below.36 in.21 in.150 degrees18 in.48 in.a. What is the area, in square inches, of the entire trapezoidal rear window? Show or explain how you got your answer.b. What fractional part of a complete circle is cleared on the rear window by the 18-inch wiper? Show or explain how you got your answer.c. What is the area, in square inches, of the part of the rear window that is cleared by the wiper? Show or explain how you got your answer.d. 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